Diffusion models, rooted in probabilistic generative modeling, are powerful tools for data generation. Initially in machine learning research, their history dates back to the mid-2010s when Denoising Autoencoders were developed. Today, they have gained prominence for their ability to generate high-quality images from text by modeling the denoising process. Current usage is in image synthesis, text generation, anomaly detection, finding utility in art, natural language processing, and cybersecurity. The future scope of diffusion models holds the potential for revolutionizing content creation, improving language understanding, making them a pivotal part of AI technologies, and solving real-world challenges. In this article, we will understand the basics of the diffusion model. Our focus will be on latent diffusion models related to text-to-image generation. We will learn to use image generation with the diffusion model in Python the Stable Diffusion model by Dream Studio. So let’s get started!
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Diffusion models belong to the class of generative models, meaning they can generate data similar to the one on which they are trained. In essence, the diffusion models destroy training data by adding noise and then learning to recover the training data by removing the noise. In the process, it learns the parameters of the neural network. We can then use this trained model and generate new data similar to training data by randomly sampling noise through the learned denoising process. This concept is similar to Variational Autoencoders (VAEs) in which we try to optimize a cost function by first projecting the data onto the latent space and then recovering it back to the starting state. In diffusion models, the system aims to model a series of noise distributions in a Markov Chain and “decodes” the data by undoing/denoising the data in a hierarchical fashion.
A diffusion denoising process modeling basically involves 2 major steps – the forward diffusion process (adding noise) and reverse diffusion process (removing noise). Let us try to understand each step one by one.
The below are the steps for forward diffusion:
Many open-source contributors collaborated to create the Stable Diffusion model, which is one of the most popular and efficient diffusion models available. It runs seamlessly on limited compute resources. It’s architecture consists of 4 components :-
1. Variational Autoencoders (VAE): Utilise it to decode pictures and translate them from latent space into pixel space. The latent space is a condensed representation of a picture that highlights its key elements. Working with latent embeddings is computationally lot cheaper and compress the latent spaces (have significantly lower dimensionality).
2. Text encoder and Tokenizer: To encode the user specific text prompt which is to generate the image.
3. The U-Net Model: Latent image representations are denoised using it. Like an autoencoder, a U-Net has a contracting path and an expanding path. A U-Net does, however, have skip connections. These aid in the information propagation from the prior layers, which helps to solve the issue of disappearing gradients. Additionally, since we ultimately lose information in the contractive path, it aids in maintaining the finer details.
In the below python implementation we will use the stable diffusion model to generate images.
!pip install transformers diffusers accelerate
!pip install xformers
from diffusers import StableDiffusionPipeline
import torch
Here we load the specific stable diffusion model in model_id below which is on Hugging face library.
model_id = "dreamlike-art/dreamlike-photoreal-2.0"
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
Here we generate 3 prompts for images we create 2 images of Alice in Wonderland with different styles and a third image of chesire cat.
prompts = ["Alice in Wonderland, Ultra HD, realistic, futuristic, detailed, octane render, photoshopped, photorealistic, soft, pastel, Aesthetic, Magical background",
"Anime style Alice in Wonderland, 90's vintage style, digital art, ultra HD, 8k, photoshopped, sharp focus, surrealism, akira style, detailed line art",
"Beautiful, abstract art of Chesire cat of Alice in wonderland, 3D, highly detailed, 8K, aesthetic"]
images = []
for i, prompt in enumerate(prompts):
image = pipe(prompt).images[0]
image.save(f'picture_{i}.jpg')
images.append(image)
In the realm of AI, researchers are currently exploring the powerful potential of diffusion models for wider application across various domains. Product designers and illustrators are experimenting with these models to quickly generate innovative prototype designs. Furthermore, several other robust models exist for generating more detailed images and can find utility in various photography tasks. Experts believe that these models will have a pivotal role in generating video content for influencers in the future.
A. There are a number of powerful diffusion models available like DALLE 2 by Open AI , Imagen by Google , Midjourney and Stable Diffusion by StabilityAI.
A. Stable Diffusion by StabilityAI is only free open source available currently.
A. There are various generative models for image generation they are GANs, VAEs, Deep Flow based models.
A. Stability AI allows user to experiment and generate images on the website by signing up on their page https://beta.dreamstudio.ai/generate . Initially, it offers free credits to its new users, and then it charges for every image generation.
A. Yes, apart from texts, we can also upload another image as a reference or edit the image by giving a prompt to remove specific objects from image or color the black and white image, etc. This service is by the RunawayML platform Image2Image
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